prompt

Optional. The number of inference steps. Defaults to 40

n_steps

Optional. The proportion of `n_steps` to run the base model for. Must be a number between 0 and 1. Defaults to 0.8

high_noise_frac


Stable Diffusion XL 1.0 base combined with the refiner as an ensemble of experts.

![](https://raw.githubusercontent.com/Stability-AI/generative-models/477d8b9a7730d9b2e92b326a770c0420d00308c9/assets/000.jpg)

Options that can be passed when loading the model:

 - `model.compile` (bool): whether or not to use `torch.compile`. Defaults to false
 - `model.enable_cpu_offload` (bool): whether or not to use [`enable_model_cpu_offload()`](https://huggingface.co/docs/diffusers/main/en/api/pipelines/stable_diffusion/stable_diffusion_xl#diffusers.StableDiffusionXLPipeline.enable_model_cpu_offload). Defaults to false


image

quickstart

Stable Diffusion XL 1.0 + refiner as an Ensemble of Experts

Text-to-Image Generation is a task in computer vision and natural language processing where the goal is to generate an image that corresponds to a given textual description.

Carton

Text-to-image generation

stabilityai / sdxl